Stable Diffusion XL – SDXL 1.0 Model
Stable Diffusion XL or SDXL is the latest image generation model that is tailored towards more photorealistic outputs with more detailed imagery and composition compared to previous SD models, including SD 2.1.
With Stable Diffusion XL you can now make more realistic images with improved face generation, produce legible text within images, and create more aesthetically pleasing art using shorter prompts.
Stable Diffusion XL Examples
SDXL 1.0 is out and it’s awesome!
Check out the best examples, or create your own!
(Both SDXL 0.9 and SDXL 1.0 examples are included)
What is Stable Diffusion XL or SDXL
Stable Diffusion XL (SDXL), is the latest AI image generation model that can generate realistic faces, legible text within the images, and better image composition, all while using shorter and simpler prompts. Just like its predecessors, SDXL has the ability to generate image variations using image-to-image prompting, inpainting (reimagining of the selected parts of an image), and outpainting (creating new parts that lie outside the image borders).
Stable Diffusion XL is currently available via DreamStudio and other image generation apps like NightCafe Studio and ClipDrop. Stable Diffusion XL model is available for download at HuggingFace.
Why is Stable Diffusion XL better
Legible text
One of the most significant differences between SDXL and the rest of AI image generation models (except for Deep Floyd), is that it can generate legible text. It wasn’t possible in the previous SD or any other models.
The text on the images is not always right, but it is lightyears ahead of most of the other AI art models.
Better human anatomy
Previous SD models had visible problems with generating correct human anatomy. They often had extra or missing limbs, extremely deformed faces, etc. Stable Diffusion XL fixes this to an extent. The quality is much better now. We hope it’ll be even better when it gets released.
Artistic styles
These don’t make SDXL better, but they are somewhat different. Stable Diffusion XL offers several artistic styles for image generation. They are: No style, Enhance, Anime, Photographic, Digital Art, Comic book, Fantasy art, Analog film, Neon punk, Isomteric, Low poly, Origami, Line Art, Craft clay, Cinematic, 3D model, and Pixel Art.
Shorter prompts
SDXL understands short prompts much better than the previous models. You don’t need lengthy walls of texts for prompts to get the desired results. Styles help achieve that to a degree, but even without them, SDXL understands you better!
Improved composition
Enhanced image composition allows for creating stunning visuals for almost any type of prompts without too much hustle.
Stable Diffusion XL API Playground
How to use the Stable Diffusion XL model
Stability AI has finally released the SDXL model on HuggingFace! You can now download the model.
Learn more about how to use the Stable Diffusion XL model offline using ComfyUI or Automatic1111.
If you don’t want to waste time installing the software and figuring out how to set it all up, you can try SDXL online on Discord or on ClipDrop, but it’ll be just a base model.
Download Stable Diffusion XL
Finally, the day has come. Stability AI has released the SDXL model into the wild. The model is available for download on HuggingFace. Click here to download the SDXL 1.0 base model.
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